huggingface / huggingface/diffusers

How to use Playground2.5 to train lora with own dataset to generate pictures of a specific style?

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bug stale
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Description

### Describe the bug

Hi,

I have been working on training models using the same dataset as "stabilityai/stable-diffusion-xl-base-1.0" with the script examples/text_to_image/train_text_to_image_lora_sdxl.py, and I achieved quite promising results.

Now, I am trying to further improve the performance by switching to Dreambooth. I am currently using playground2.5 with examples/dreambooth/train_dreambooth_lora_sdxl.py. However, after multiple parameter tuning attempts, the performance is still not as good as the SDXL base model.

I am unsure what might be causing this.

### Reproduction

![image](https://github.com/user-attachments/assets/339a0e9b-de08-408d-a43a-495f86b5e1df)

### Logs

_No response_

### System Info

- 🤗 Diffusers version: 0.31.0.dev0
- Platform: Linux-5.14.0-284.25.1.el9_2.x86_64-x86_64-with-glibc2.17
- Running on Google Colab?: No
- Python version: 3.8.20
- PyTorch version (GPU?): 2.2.0 (True)
- Flax version (CPU?/GPU?/TPU?): not installed (NA)
- Jax version: not installed
- JaxLib version: not installed
- Huggingface_hub version: 0.25.2
- Transformers version: 4.45.2
- Accelerate version: 1.0.1
- PEFT version: 0.13.2
- Bitsandbytes version: 0.44.1
- Safetensors version: 0.4.5
- xFormers version: not installed
- Accelerator: NVIDIA H800, 81559 MiB
- Using GPU in script?:
- Using distributed or parallel set-up in script?:

### Who can help?

_No response_

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